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In this paper, we study the existence of infinitely many periodic solutions for the non-autonomous second-order Hamiltonian systems with symmetry. Based on the minimax methods in critical point theory, in particular, the fountain theorem of Bartsch and the symmetric mountain pass lemma due to Kajikiya, we obtain the existence results for both the superquadratic case and the subquadratic case, which unify and sharply improve some recent results in the literature.
In this paper, the bond-based peridynamic system is analysed as a non-local boundary-value problem with volume constraint. The study extends earlier works in the literature on non-local diffusion and non-local peridynamic models, to include non-positive definite kernels. We prove the well-posedness of both linear and nonlinear variational problems with volume constraints. The analysis is based on some non-local Poincaré-type inequalities and the compactness of the associated non-local operators. It also offers careful characterizations of the associated solution spaces, such as compact embedding, separability and completeness. In the limit of vanishing non-locality, the convergence of the peridynamic system to the classical Navier equations of elasticity with Poisson ratio ¼ is demonstrated.
Various modes of convergence of measurable functions to a limit function were considered in Chapter 6, and will be restated here with the special terminology customarily used in the probabilistic context. In this section the modes of convergence all concern a sequence {ξn} of r.v.'s on the same probability space (Ω, F, P) such that the values ξn(ω) “become close” (in some “local” or “global” sense) to a “limiting r.v.” ξ(ω)as n → ∞. In the next section we shall consider the weaker form of convergence where the ξn's can be defined on different spaces, and where one is interested in only the limiting form of the distribution of the ξn (i.e. Pξ-1n B for Borel sets B). This “convergence in distribution” has wide use in statistical theory and application.
The later sections of the chapter will be concerned with various important relationships between the forms of convergence, convergence of series of independent r.v.'s, and related topics. Note that in certain calculations concerning convergence (especially in Section 11.5) it will be implicitly assumed that the r.v.'s involved are defined for all ω. No comment will be made in these cases, since it is a trivial matter to obtain these results for r.v.'s ξn not defined everywhere by considering ξ*n defined for all ω, and equal to ξna.s.
Relaxation of the requirement of a measure that it be nonnegative yields what is usually called a signed measure. Specifically this is an extended real-valued, countably additive set function μ on a class ε (containing ∅), such that μ(∅) = 0, and such that μ assumes at most one of the values +∞ and −∞ on ε. As for measures, a signed measure μ defined on a class ε, is called finite on ε if |μ(E)| < ∞, for each E ∈ ε, and σ-finite if for each E ∈ ε there is a sequence {En}∞n=1 of sets in ε with E ⊂ ∪∞n=1En and |μ(En)| < ∞, that is, if E can be covered by the union of a sequence of sets with finite (signed) measure. It will usually be assumed that the class on which μ is defined is a σ-ring or σ-field.
Some of the important properties of measures (see Section 2.2) hold also for signed measures. In particular a signed measure is subtractive and continuous from below and above. The basic properties of signed measures are given in the following theorem.
Theorem 5.1.1Let μ be a signed measure on a σ-ring S.
We shall consider sets consisting of elements or points. The nature of the points will be left unspecified – examples are points in a Euclidean space, sequences of numbers, functions, elementary events, etc. Small letters will be used for points.
Sets are aggregates or collections of such points. Capital letters will be used for sets.
A set is defined by a property. That is, given a point, there is a criterion to decide whether it belongs to a given set, e.g. the set which is the open interval (−1, 1) on the real line is defined by the property that it contains a point x if and only if |x| < 1.
A set may be written as {x: P(x)} where P(x) is the property defining the set; e.g. {x: |x| < 1} is the above set consisting of all points x for which |x| < 1, i.e. (−1, 1).
In any given situation, all the points considered will belong to a fixed set called the whole space and usually denoted by X. This assumption avoids some difficulties which arise in the logical foundations of set theory.
Classes or collections of sets are just aggregates whose elements themselves are sets, e.g. the class of all intervals of the real line, the class of all circles in the plane whose centers are at the origin, and so on. Script capitals will be used for classes of sets.
Our aim in this final chapter is to indicate how basic distributional theory for stochastic processes, alias random functions, may be developed from the considerations of Chapters 7 and 9. This is primarily for reference and for readers with a potential interest in the topic. The theory will be first illustrated by a discussion of the definition of the Wiener process, and conditions for sample function continuity. This will be complemented, and the chapter completed with a sketch of construction and basic properties of point processes and random measures in a purely measure-theoretic framework, consistent with the nontopological flavor of the entire volume.
Random functions and stochastic processes
In this section we introduce some basic distributional theory for stochastic processes and random functions, using the product space measures of Chapter 7 and the random element concepts of Chapter 9.
By a stochastic process one traditionally means a family of real random variables {ξt: t ∈ T} (ξ, = ξt(ω)) on a probability space (Ω, F, P), T being a set indexing the ξt. If T = {1,2,3,…} or {…, −2, −1,0,1,2,…} the family {ξn: n =1,2,…} or {ξn: n = …,−2, −1,0,1,2,∧…} is referred to as a stochastic sequence or discrete parameter stochastic process, whereas {ξt: t ∈ T} is termed a continuous parameter stochastic process if T is an interval (finite or infinite).
James constructed a functorial homotopy decomposition for path-connected, p ointed CW-complexes X. We generalize this to a p-local functorial decomposition of ΣA, where A is any functorial retract of a looped co-H-space. This is used to construct Hopf invariants in a more general context. In addition, when A = ΩY is the loops space of a co-H-space, we show that the wedge summands of ΣΩY further functorially decompose by using an action of an appropriate symmetric group. As a valuable example, we give an application to the theory of quasi-symmetric functions.
We consider a family of bounded dissipative asymptotically compact semigroups depending on a parameter, and study the continuity properties of the corresponding family of its global attractors. We exploit the idea of the uniform exponential attraction property to discuss the continuity properties of the family of attractors and estimate the rate of convergence of the approximating attractors to the limit one. Showing a range of applications of an abstract framework, we focus much of our attention on a perturbed damped wave equation. In this latter case our results involve nonlinearities with critical exponents, for which the continuity of the family of attractors is concluded, including the rate of convergence and the regularity of the limit attractor. This complements the results in the literature.
This work arises from lecture notes for a two semester basic course sequence in Measure and Probability Theory given for first year Statistics graduate students at the University of North Carolina, evolving through many generations of handwritten, typed, mimeographed, and finally LaTeX editions. Their focus is to provide basic course material, tailored to the background of our students, and influenced very much by their reactions and the changing emphases of the years. We see this as one side of an avowed department educational mission to provide solid and diverse basic course training common to all our students, who will later specialize in diverse areas from the very theoretical to the very applied.
The notes originated in the 1960's from a “Halmos style” measure theory course. As may be apparent (to those of sufficient age) the measure theory section has preserved that basic flavor with numerous obvious modernizations (beginning with the early use of the Sierpinski-type classes more suited than monotone class theorems for probabilistic applications), and exposition more tailored to the particular audience. Even the early “Halmos framework” of rings and σ-rings has been retained up to a point since these notions are useful in applications (e.g. point process theory) and their inclusion requires no significant further effort. Integration itself is discussed within the customary σ-field framework so the students have no difficulty in relating to other works.
Employing the Kolodner–Coffman method, we show the exact multiplicity of positive solutions for the one-dimensional p-Laplacian that is subject to a Dirichlet boundary condition with a positive convex nonlinearity and an indefinite weight function.
A well-known theorem proved by Lazer and Solimini claims that the singular equation
has a periodic solution if and only if the mean value of the continuous external force is positive. In this paper, we show that this result cannot be extended to the case when h is an integrable function, unless additional assumptions are introduced. In addition, for each p ≥ 1 and h-integrable function in the pth power, we give a sharp condition guaranteeing the existence of periodic solutions to the above-mentioned equation, showing that there is a close relation between p and the order of the singularity λ.
In this paper, the existence of at least three non-negative solutions to non-local boundary-value problems for second-order differential equations with deviating arguments α and ϛ is investigated. Sufficient conditions, which guarantee the existence of positive solutions, are obtained using the Avery–Peterson theorem. We discuss our problem for both advanced and delayed arguments. An example is added to illustrate the results.
With financial modelling requiring a better understanding of model risk, it is helpful to be able to vary assumptions about underlying probability distributions in an efficient manner, preferably without the noise induced by resampling distributions managed by Monte Carlo methods. This paper presents differential equations and solution methods for the functions of the form Q(x) = F−1(G(x)), where F and G are cumulative distribution functions. Such functions allow the direct recycling of Monte Carlo samples from one distribution into samples from another. The method may be developed analytically for certain special cases, and illuminate the idea that it is a more precise form of the traditional Cornish–Fisher expansion. In this manner the model risk of distributional risk may be assessed free of the Monte Carlo noise associated with resampling. The method may also be regarded as providing both analytical and numerical bases for doing more precise Cornish–Fisher transformations. Examples are given of equations for converting normal samples to Student t, and converting exponential to normal. In the case of the normal distribution, the change of variables employed allows the sampling to take place to good accuracy based on a single rational approximation over a very wide range of sample space. The avoidance of branching statements is of use in optimal graphics processing unit (GPU) computations as it avoids the effect of branch divergence. We give a branch-free normal quantile that offers performance improvements in a GPU environment while retaining the best precision characteristics of well-known methods. We also offer models with low probability branch divergence. Comparisons of new and existing forms are made on Nvidia GeForce GTX Titan and Tesla C2050 GPUs. We argue that in both single- and double-precisions, the change-of-variables approach offers the most GPU-optimal Gaussian quantile yet, working faster than the Cuda 5.5 built-in function.
In a 1976 paper published in Science, Knuth presented an algorithm to sample (non-uniform) self-avoiding walks crossing a square of side k. From this sample, he constructed an estimator for the number of such walks. The quality of this estimator is directly related to the (relative) variance of a certain random variable Xk. From his experiments, Knuth suspected that this variance was extremely large (so that the estimator would not be very efficient). But how large? For the analogous Rosenbluth algorithm, which samples unconfined self-avoiding walks of length n, the variance of the corresponding estimator is believed to be exponential in n.
A few years ago, Bassetti and Diaconis showed that, for a sampler à la Knuth that generates walks crossing a k × k square and consisting of North and East steps, the relative variance is only $O(\sqrt k)$. In this note we take one step further and show that, for walks consisting of North, South and East steps, the relative variance jumps to $2^{k(k+1)}/(k+1)^{2k}$. This is exponential in the average length of the walks, which is of order k2. We also obtain partial results for general self-avoiding walks crossing a square, suggesting that the relative variance could be exponential in k2 (which is again the average length of these walks).
Knuth's algorithm is a basic example of a widely used technique called sequential importance sampling. The present paper, following the paper by Bassetti and Diaconis, is one of very few examples where the variance of the estimator can be found.
In this work we investigate limiting values of the lift and drag coefficients of profiles in the Helmholtz–Kirchhoff (infinite cavity) flow. The coefficients are based on the wetted arc length of profile surfaces. The problem is to find global minimum and maximum values of the drag coefficient CD under a given lift coefficient CL. We reduce the problem to a constrained problem of calculus of variations and solve it analytically. In so doing we do not only determine extremals but also strictly prove that these extremals realize global extrema. The proofs are based on non-trivial application of Jensen's inequality. The solution of the problem allows us to construct the domain of possible variations of coefficients CL and CD and define maximum and minimum values of the lift-to-drag ratios CL/CD for a given CL.
In the previous chapter we saw that every (closed) invariant subspace of Dμ is cyclic, in other words, that it is generated by a single function in Dμ. In this chapter, we shall take the opposite point of view: starting with f ∈ Dμ, can we identify the invariant subspace that it generates? It is easy to see that g belongs to this invariant subspace if and only if there is a sequence of polynomials (pn) such that ||pnf − g||Dμ → 0, but in practice it is often difficult to determine which g have this property. Therefore we seek other, more explicit descriptions of the invariant subspace generated by f.
In particular, we pose the question: which functions f generate the whole of Dμ? A complete answer to this is not known, even for the special case of the Dirichlet space D. In fact the characterization of those functions cyclic for D is the subject of a conjecture involving the logarithmic capacity of boundary zero sets. We shall present some partial solutions to this conjecture.
Cyclicity inDμ
We begin by formalizing the notions above. Let μ be a finite positive measure on T, and let Dμ be the associated harmonically weighted Dirichlet space. As usual, if μ is normalized Lebesgue measure, then Dμ is just the classical Dirichlet space D, so all the results of this section apply in particular to D.